1350 lines
53 KiB
ReStructuredText
1350 lines
53 KiB
ReStructuredText
Text-to-Image Generation with Stable Diffusion v2 and OpenVINO™
|
||
===============================================================
|
||
|
||
Stable Diffusion v2 is the next generation of Stable Diffusion model a
|
||
Text-to-Image latent diffusion model created by the researchers and
|
||
engineers from `Stability AI <https://stability.ai/>`__ and
|
||
`LAION <https://laion.ai/>`__.
|
||
|
||
General diffusion models are machine learning systems that are trained
|
||
to denoise random gaussian noise step by step, to get to a sample of
|
||
interest, such as an image. Diffusion models have shown to achieve
|
||
state-of-the-art results for generating image data. But one downside of
|
||
diffusion models is that the reverse denoising process is slow. In
|
||
addition, these models consume a lot of memory because they operate in
|
||
pixel space, which becomes unreasonably expensive when generating
|
||
high-resolution images. Therefore, it is challenging to train these
|
||
models and also use them for inference. OpenVINO brings capabilities to
|
||
run model inference on Intel hardware and opens the door to the
|
||
fantastic world of diffusion models for everyone!
|
||
|
||
In previous notebooks, we already discussed how to run `Text-to-Image
|
||
generation and Image-to-Image generation using Stable Diffusion
|
||
v1 <stable-diffusion-text-to-image-with-output.html>`__
|
||
and `controlling its generation process using
|
||
ControlNet <./controlnet-stable-diffusion/controlnet-stable-diffusion.ipynb>`__.
|
||
Now is turn of Stable Diffusion v2.
|
||
|
||
Stable Diffusion v2: What’s new?
|
||
--------------------------------
|
||
|
||
The new stable diffusion model offers a bunch of new features inspired
|
||
by the other models that have emerged since the introduction of the
|
||
first iteration. Some of the features that can be found in the new model
|
||
are:
|
||
|
||
- The model comes with a new robust encoder, OpenCLIP, created by LAION
|
||
and aided by Stability AI; this version v2 significantly enhances the
|
||
produced photos over the V1 versions.
|
||
- The model can now generate images in a 768x768 resolution, offering
|
||
more information to be shown in the generated images.
|
||
- The model finetuned with
|
||
`v-objective <https://arxiv.org/abs/2202.00512>`__. The
|
||
v-parameterization is particularly useful for numerical stability
|
||
throughout the diffusion process to enable progressive distillation
|
||
for models. For models that operate at higher resolution, it is also
|
||
discovered that the v-parameterization avoids color shifting
|
||
artifacts that are known to affect high resolution diffusion models,
|
||
and in the video setting it avoids temporal color shifting that
|
||
sometimes appears with epsilon-prediction used in Stable Diffusion
|
||
v1.
|
||
- The model also comes with a new diffusion model capable of running
|
||
upscaling on the images generated. Upscaled images can be adjusted up
|
||
to 4 times the original image. Provided as separated model, for more
|
||
details please check
|
||
`stable-diffusion-x4-upscaler <https://huggingface.co/stabilityai/stable-diffusion-x4-upscaler>`__
|
||
- The model comes with a new refined depth architecture capable of
|
||
preserving context from prior generation layers in an image-to-image
|
||
setting. This structure preservation helps generate images that
|
||
preserving forms and shadow of objects, but with different content.
|
||
- The model comes with an updated inpainting module built upon the
|
||
previous model. This text-guided inpainting makes switching out parts
|
||
in the image easier than before.
|
||
|
||
This notebook demonstrates how to convert and run Stable Diffusion v2
|
||
model using OpenVINO.
|
||
|
||
Notebook contains the following steps:
|
||
|
||
1. Create PyTorch models pipeline using Diffusers library.
|
||
2. Convert PyTorch models to OpenVINO IR format, using model conversion
|
||
API.
|
||
3. Apply hybrid post-training quantization to UNet model with
|
||
`NNCF <https://github.com/openvinotoolkit/nncf/>`__.
|
||
4. Run Stable Diffusion v2 Text-to-Image pipeline with OpenVINO.
|
||
|
||
**Note:** This is the full version of the Stable Diffusion text-to-image
|
||
implementation. If you would like to get started and run the notebook
|
||
quickly, check out `stable-diffusion-v2-text-to-image-demo
|
||
notebook <stable-diffusion-v2-with-output.html>`__.
|
||
|
||
Table of contents:
|
||
^^^^^^^^^^^^^^^^^^
|
||
|
||
- `Prerequisites <#prerequisites>`__
|
||
- `Stable Diffusion v2 for Text-to-Image
|
||
Generation <#stable-diffusion-v2-for-text-to-image-generation>`__
|
||
|
||
- `Stable Diffusion in Diffusers
|
||
library <#stable-diffusion-in-diffusers-library>`__
|
||
- `Convert models to OpenVINO Intermediate representation (IR)
|
||
format <#convert-models-to-openvino-intermediate-representation-ir-format>`__
|
||
- `Text Encoder <#text-encoder>`__
|
||
- `U-Net <#u-net>`__
|
||
- `VAE <#vae>`__
|
||
- `Prepare Inference Pipeline <#prepare-inference-pipeline>`__
|
||
- `Configure Inference Pipeline <#configure-inference-pipeline>`__
|
||
|
||
- `Quantization <#quantization>`__
|
||
|
||
- `Prepare calibration dataset <#prepare-calibration-dataset>`__
|
||
- `Run Hybrid Model Quantization <#run-hybrid-model-quantization>`__
|
||
- `Compare inference time of the FP16 and INT8
|
||
pipelines <#compare-inference-time-of-the-fp16-and-int8-pipelines>`__
|
||
|
||
- `Run Text-to-Image generation <#run-text-to-image-generation>`__
|
||
|
||
Prerequisites
|
||
-------------
|
||
|
||
|
||
|
||
install required packages
|
||
|
||
.. code:: ipython3
|
||
|
||
%pip install -q "diffusers>=0.14.0" "openvino>=2023.1.0" "datasets>=2.14.6" "transformers>=4.25.1" "gradio>=4.19" "torch>=2.1" Pillow opencv-python --extra-index-url https://download.pytorch.org/whl/cpu
|
||
%pip install -q "nncf>=2.9.0"
|
||
|
||
|
||
.. parsed-literal::
|
||
|
||
Note: you may need to restart the kernel to use updated packages.
|
||
|
||
|
||
Stable Diffusion v2 for Text-to-Image Generation
|
||
------------------------------------------------
|
||
|
||
|
||
|
||
To start, let’s look on Text-to-Image process for Stable Diffusion v2.
|
||
We will use `Stable Diffusion
|
||
v2-1 <https://huggingface.co/stabilityai/stable-diffusion-2-1>`__ model
|
||
for these purposes. The main difference from Stable Diffusion v2 and
|
||
Stable Diffusion v2.1 is usage of more data, more training, and less
|
||
restrictive filtering of the dataset, that gives promising results for
|
||
selecting wide range of input text prompts. More details about model can
|
||
be found in `Stability AI blog
|
||
post <https://stability.ai/blog/stablediffusion2-1-release7-dec-2022>`__
|
||
and original model
|
||
`repository <https://github.com/Stability-AI/stablediffusion>`__.
|
||
|
||
Stable Diffusion in Diffusers library
|
||
~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~
|
||
|
||
To work with Stable Diffusion
|
||
v2, we will use Hugging Face
|
||
`Diffusers <https://github.com/huggingface/diffusers>`__ library. To
|
||
experiment with Stable Diffusion models, Diffusers exposes the
|
||
`StableDiffusionPipeline <https://huggingface.co/docs/diffusers/using-diffusers/conditional_image_generation>`__
|
||
similar to the `other Diffusers
|
||
pipelines <https://huggingface.co/docs/diffusers/api/pipelines/overview>`__.
|
||
The code below demonstrates how to create ``StableDiffusionPipeline``
|
||
using ``stable-diffusion-2-1``:
|
||
|
||
.. code:: ipython3
|
||
|
||
from diffusers import StableDiffusionPipeline
|
||
|
||
pipe = StableDiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2-1-base").to("cpu")
|
||
|
||
# for reducing memory consumption get all components from pipeline independently
|
||
text_encoder = pipe.text_encoder
|
||
text_encoder.eval()
|
||
unet = pipe.unet
|
||
unet.eval()
|
||
vae = pipe.vae
|
||
vae.eval()
|
||
|
||
conf = pipe.scheduler.config
|
||
|
||
del pipe
|
||
|
||
|
||
|
||
.. parsed-literal::
|
||
|
||
Loading pipeline components...: 0%| | 0/6 [00:00<?, ?it/s]
|
||
|
||
|
||
Convert models to OpenVINO Intermediate representation (IR) format
|
||
~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~
|
||
|
||
|
||
|
||
Starting from 2023.0 release, OpenVINO supports PyTorch models directly
|
||
via Model Conversion API. ``ov.convert_model`` function accepts instance
|
||
of PyTorch model and example inputs for tracing and returns object of
|
||
``ov.Model`` class, ready to use or save on disk using ``ov.save_model``
|
||
function.
|
||
|
||
The pipeline consists of three important parts:
|
||
|
||
- Text Encoder to create condition to generate an image from a text
|
||
prompt.
|
||
- U-Net for step-by-step denoising latent image representation.
|
||
- Autoencoder (VAE) for decoding latent space to image.
|
||
|
||
Let us convert each part:
|
||
|
||
Text Encoder
|
||
~~~~~~~~~~~~
|
||
|
||
|
||
|
||
The text-encoder is responsible for transforming the input prompt, for
|
||
example, “a photo of an astronaut riding a horse” into an embedding
|
||
space that can be understood by the U-Net. It is usually a simple
|
||
transformer-based encoder that maps a sequence of input tokens to a
|
||
sequence of latent text embeddings.
|
||
|
||
The input of the text encoder is tensor ``input_ids``, which contains
|
||
indexes of tokens from text processed by the tokenizer and padded to the
|
||
maximum length accepted by the model. Model outputs are two tensors:
|
||
``last_hidden_state`` - hidden state from the last MultiHeadAttention
|
||
layer in the model and ``pooler_out`` - pooled output for whole model
|
||
hidden states.
|
||
|
||
.. code:: ipython3
|
||
|
||
from pathlib import Path
|
||
|
||
sd2_1_model_dir = Path("sd2.1")
|
||
sd2_1_model_dir.mkdir(exist_ok=True)
|
||
|
||
.. code:: ipython3
|
||
|
||
import gc
|
||
import torch
|
||
import openvino as ov
|
||
|
||
TEXT_ENCODER_OV_PATH = sd2_1_model_dir / "text_encoder.xml"
|
||
|
||
|
||
def cleanup_torchscript_cache():
|
||
"""
|
||
Helper for removing cached model representation
|
||
"""
|
||
torch._C._jit_clear_class_registry()
|
||
torch.jit._recursive.concrete_type_store = torch.jit._recursive.ConcreteTypeStore()
|
||
torch.jit._state._clear_class_state()
|
||
|
||
|
||
def convert_encoder(text_encoder: torch.nn.Module, ir_path: Path):
|
||
"""
|
||
Convert Text Encoder model to IR.
|
||
Function accepts pipeline, prepares example inputs for conversion
|
||
Parameters:
|
||
text_encoder (torch.nn.Module): text encoder PyTorch model
|
||
ir_path (Path): File for storing model
|
||
Returns:
|
||
None
|
||
"""
|
||
if not ir_path.exists():
|
||
input_ids = torch.ones((1, 77), dtype=torch.long)
|
||
# switch model to inference mode
|
||
text_encoder.eval()
|
||
|
||
# disable gradients calculation for reducing memory consumption
|
||
with torch.no_grad():
|
||
# export model
|
||
ov_model = ov.convert_model(
|
||
text_encoder, # model instance
|
||
example_input=input_ids, # example inputs for model tracing
|
||
input=([1, 77],), # input shape for conversion
|
||
)
|
||
ov.save_model(ov_model, ir_path)
|
||
del ov_model
|
||
cleanup_torchscript_cache()
|
||
print("Text Encoder successfully converted to IR")
|
||
|
||
|
||
if not TEXT_ENCODER_OV_PATH.exists():
|
||
convert_encoder(text_encoder, TEXT_ENCODER_OV_PATH)
|
||
else:
|
||
print(f"Text encoder will be loaded from {TEXT_ENCODER_OV_PATH}")
|
||
|
||
del text_encoder
|
||
gc.collect();
|
||
|
||
|
||
.. parsed-literal::
|
||
|
||
Text encoder will be loaded from sd2.1/text_encoder.xml
|
||
|
||
|
||
U-Net
|
||
~~~~~
|
||
|
||
|
||
|
||
U-Net model gradually denoises latent image representation guided by
|
||
text encoder hidden state.
|
||
|
||
U-Net model has three inputs:
|
||
|
||
- ``sample`` - latent image sample from previous step. Generation
|
||
process has not been started yet, so you will use random noise.
|
||
- ``timestep`` - current scheduler step.
|
||
- ``encoder_hidden_state`` - hidden state of text encoder.
|
||
|
||
Model predicts the ``sample`` state for the next step.
|
||
|
||
Generally, U-Net model conversion process remain the same like in Stable
|
||
Diffusion v1, expect small changes in input sample size. Our model was
|
||
pretrained to generate images with resolution 768x768, initial latent
|
||
sample size for this case is 96x96. Besides that, for different use
|
||
cases like inpainting and depth to image generation model also can
|
||
accept additional image information: depth map or mask as channel-wise
|
||
concatenation with initial latent sample. For converting U-Net model for
|
||
such use cases required to modify number of input channels.
|
||
|
||
.. code:: ipython3
|
||
|
||
import numpy as np
|
||
|
||
UNET_OV_PATH = sd2_1_model_dir / "unet.xml"
|
||
|
||
|
||
def convert_unet(
|
||
unet: torch.nn.Module,
|
||
ir_path: Path,
|
||
num_channels: int = 4,
|
||
width: int = 64,
|
||
height: int = 64,
|
||
):
|
||
"""
|
||
Convert Unet model to IR format.
|
||
Function accepts pipeline, prepares example inputs for conversion
|
||
Parameters:
|
||
unet (torch.nn.Module): UNet PyTorch model
|
||
ir_path (Path): File for storing model
|
||
num_channels (int, optional, 4): number of input channels
|
||
width (int, optional, 64): input width
|
||
height (int, optional, 64): input height
|
||
Returns:
|
||
None
|
||
"""
|
||
dtype_mapping = {torch.float32: ov.Type.f32, torch.float64: ov.Type.f64}
|
||
if not ir_path.exists():
|
||
# prepare inputs
|
||
encoder_hidden_state = torch.ones((2, 77, 1024))
|
||
latents_shape = (2, num_channels, width, height)
|
||
latents = torch.randn(latents_shape)
|
||
t = torch.from_numpy(np.array(1, dtype=np.float32))
|
||
unet.eval()
|
||
dummy_inputs = (latents, t, encoder_hidden_state)
|
||
input_info = []
|
||
for input_tensor in dummy_inputs:
|
||
shape = ov.PartialShape(tuple(input_tensor.shape))
|
||
element_type = dtype_mapping[input_tensor.dtype]
|
||
input_info.append((shape, element_type))
|
||
|
||
with torch.no_grad():
|
||
ov_model = ov.convert_model(unet, example_input=dummy_inputs, input=input_info)
|
||
ov.save_model(ov_model, ir_path)
|
||
del ov_model
|
||
cleanup_torchscript_cache()
|
||
print("U-Net successfully converted to IR")
|
||
|
||
|
||
if not UNET_OV_PATH.exists():
|
||
convert_unet(unet, UNET_OV_PATH, width=96, height=96)
|
||
del unet
|
||
gc.collect()
|
||
else:
|
||
del unet
|
||
gc.collect();
|
||
|
||
VAE
|
||
~~~
|
||
|
||
|
||
|
||
The VAE model has two parts, an encoder and a decoder. The encoder is
|
||
used to convert the image into a low dimensional latent representation,
|
||
which will serve as the input to the U-Net model. The decoder,
|
||
conversely, transforms the latent representation back into an image.
|
||
|
||
During latent diffusion training, the encoder is used to get the latent
|
||
representations (latents) of the images for the forward diffusion
|
||
process, which applies more and more noise at each step. During
|
||
inference, the denoised latents generated by the reverse diffusion
|
||
process are converted back into images using the VAE decoder. When you
|
||
run inference for Text-to-Image, there is no initial image as a starting
|
||
point. You can skip this step and directly generate initial random
|
||
noise.
|
||
|
||
When running Text-to-Image pipeline, we will see that we **only need the
|
||
VAE decoder**, but preserve VAE encoder conversion, it will be useful in
|
||
next chapter of our tutorial.
|
||
|
||
Note: This process will take a few minutes and use significant amount of
|
||
RAM (recommended at least 32GB).
|
||
|
||
.. code:: ipython3
|
||
|
||
VAE_ENCODER_OV_PATH = sd2_1_model_dir / "vae_encoder.xml"
|
||
|
||
|
||
def convert_vae_encoder(vae: torch.nn.Module, ir_path: Path, width: int = 512, height: int = 512):
|
||
"""
|
||
Convert VAE model to IR format.
|
||
VAE model, creates wrapper class for export only necessary for inference part,
|
||
prepares example inputs for onversion
|
||
Parameters:
|
||
vae (torch.nn.Module): VAE PyTorch model
|
||
ir_path (Path): File for storing model
|
||
width (int, optional, 512): input width
|
||
height (int, optional, 512): input height
|
||
Returns:
|
||
None
|
||
"""
|
||
|
||
class VAEEncoderWrapper(torch.nn.Module):
|
||
def __init__(self, vae):
|
||
super().__init__()
|
||
self.vae = vae
|
||
|
||
def forward(self, image):
|
||
return self.vae.encode(x=image)["latent_dist"].sample()
|
||
|
||
if not ir_path.exists():
|
||
vae_encoder = VAEEncoderWrapper(vae)
|
||
vae_encoder.eval()
|
||
image = torch.zeros((1, 3, width, height))
|
||
with torch.no_grad():
|
||
ov_model = ov.convert_model(vae_encoder, example_input=image, input=([1, 3, width, height],))
|
||
ov.save_model(ov_model, ir_path)
|
||
del ov_model
|
||
cleanup_torchscript_cache()
|
||
print("VAE encoder successfully converted to IR")
|
||
|
||
|
||
def convert_vae_decoder(vae: torch.nn.Module, ir_path: Path, width: int = 64, height: int = 64):
|
||
"""
|
||
Convert VAE decoder model to IR format.
|
||
Function accepts VAE model, creates wrapper class for export only necessary for inference part,
|
||
prepares example inputs for conversion
|
||
Parameters:
|
||
vae (torch.nn.Module): VAE model
|
||
ir_path (Path): File for storing model
|
||
width (int, optional, 64): input width
|
||
height (int, optional, 64): input height
|
||
Returns:
|
||
None
|
||
"""
|
||
|
||
class VAEDecoderWrapper(torch.nn.Module):
|
||
def __init__(self, vae):
|
||
super().__init__()
|
||
self.vae = vae
|
||
|
||
def forward(self, latents):
|
||
return self.vae.decode(latents)
|
||
|
||
if not ir_path.exists():
|
||
vae_decoder = VAEDecoderWrapper(vae)
|
||
latents = torch.zeros((1, 4, width, height))
|
||
|
||
vae_decoder.eval()
|
||
with torch.no_grad():
|
||
ov_model = ov.convert_model(vae_decoder, example_input=latents, input=([1, 4, width, height],))
|
||
ov.save_model(ov_model, ir_path)
|
||
del ov_model
|
||
cleanup_torchscript_cache()
|
||
print("VAE decoder successfully converted to IR")
|
||
|
||
|
||
if not VAE_ENCODER_OV_PATH.exists():
|
||
convert_vae_encoder(vae, VAE_ENCODER_OV_PATH, 768, 768)
|
||
else:
|
||
print(f"VAE encoder will be loaded from {VAE_ENCODER_OV_PATH}")
|
||
|
||
VAE_DECODER_OV_PATH = sd2_1_model_dir / "vae_decoder.xml"
|
||
|
||
if not VAE_DECODER_OV_PATH.exists():
|
||
convert_vae_decoder(vae, VAE_DECODER_OV_PATH, 96, 96)
|
||
else:
|
||
print(f"VAE decoder will be loaded from {VAE_DECODER_OV_PATH}")
|
||
|
||
del vae
|
||
gc.collect();
|
||
|
||
|
||
.. parsed-literal::
|
||
|
||
VAE encoder will be loaded from sd2.1/vae_encoder.xml
|
||
VAE decoder will be loaded from sd2.1/vae_decoder.xml
|
||
|
||
|
||
Prepare Inference Pipeline
|
||
~~~~~~~~~~~~~~~~~~~~~~~~~~
|
||
|
||
|
||
|
||
Putting it all together, let us now take a closer look at how the model
|
||
works in inference by illustrating the logical flow.
|
||
|
||
.. figure:: https://github.com/openvinotoolkit/openvino_notebooks/assets/22090501/ec454103-0d28-48e3-a18e-b55da3fab381
|
||
:alt: text2img-stable-diffusion v2
|
||
|
||
text2img-stable-diffusion v2
|
||
|
||
The stable diffusion model takes both a latent seed and a text prompt as
|
||
input. The latent seed is then used to generate random latent image
|
||
representations of size :math:`96 \times 96` where as the text prompt is
|
||
transformed to text embeddings of size :math:`77 \times 1024` via
|
||
OpenCLIP’s text encoder.
|
||
|
||
Next, the U-Net iteratively *denoises* the random latent image
|
||
representations while being conditioned on the text embeddings. The
|
||
output of the U-Net, being the noise residual, is used to compute a
|
||
denoised latent image representation via a scheduler algorithm. Many
|
||
different scheduler algorithms can be used for this computation, each
|
||
having its pros and cons. For Stable Diffusion, it is recommended to use
|
||
one of:
|
||
|
||
- `PNDM
|
||
scheduler <https://github.com/huggingface/diffusers/blob/main/src/diffusers/schedulers/scheduling_pndm.py>`__
|
||
- `DDIM
|
||
scheduler <https://github.com/huggingface/diffusers/blob/main/src/diffusers/schedulers/scheduling_ddim.py>`__
|
||
- `K-LMS
|
||
scheduler <https://github.com/huggingface/diffusers/blob/main/src/diffusers/schedulers/scheduling_lms_discrete.py>`__
|
||
|
||
Theory on how the scheduler algorithm function works is out of scope for
|
||
this notebook, but in short, you should remember that they compute the
|
||
predicted denoised image representation from the previous noise
|
||
representation and the predicted noise residual. For more information,
|
||
it is recommended to look into `Elucidating the Design Space of
|
||
Diffusion-Based Generative Models <https://arxiv.org/abs/2206.00364>`__.
|
||
|
||
The chart above looks very similar to Stable Diffusion V1 from
|
||
`notebook <stable-diffusion-text-to-image-with-output.html>`__,
|
||
but there is some small difference in details:
|
||
|
||
- Changed input resolution for U-Net model.
|
||
- Changed text encoder and as the result size of its hidden state
|
||
embeddings.
|
||
- Additionally, to improve image generation quality authors introduced
|
||
negative prompting. Technically, positive prompt steers the diffusion
|
||
toward the images associated with it, while negative prompt steers
|
||
the diffusion away from it.In other words, negative prompt declares
|
||
undesired concepts for generation image, e.g. if we want to have
|
||
colorful and bright image, gray scale image will be result which we
|
||
want to avoid, in this case gray scale can be treated as negative
|
||
prompt. The positive and negative prompt are in equal footing. You
|
||
can always use one with or without the other. More explanation of how
|
||
it works can be found in this
|
||
`article <https://stable-diffusion-art.com/how-negative-prompt-work/>`__.
|
||
|
||
.. code:: ipython3
|
||
|
||
import inspect
|
||
from typing import List, Optional, Union, Dict
|
||
|
||
import PIL
|
||
import cv2
|
||
import torch
|
||
|
||
from transformers import CLIPTokenizer
|
||
from diffusers import DiffusionPipeline
|
||
from diffusers.schedulers import DDIMScheduler, LMSDiscreteScheduler, PNDMScheduler
|
||
|
||
|
||
def scale_fit_to_window(dst_width: int, dst_height: int, image_width: int, image_height: int):
|
||
"""
|
||
Preprocessing helper function for calculating image size for resize with peserving original aspect ratio
|
||
and fitting image to specific window size
|
||
|
||
Parameters:
|
||
dst_width (int): destination window width
|
||
dst_height (int): destination window height
|
||
image_width (int): source image width
|
||
image_height (int): source image height
|
||
Returns:
|
||
result_width (int): calculated width for resize
|
||
result_height (int): calculated height for resize
|
||
"""
|
||
im_scale = min(dst_height / image_height, dst_width / image_width)
|
||
return int(im_scale * image_width), int(im_scale * image_height)
|
||
|
||
|
||
def preprocess(image: PIL.Image.Image):
|
||
"""
|
||
Image preprocessing function. Takes image in PIL.Image format, resizes it to keep aspect ration and fits to model input window 512x512,
|
||
then converts it to np.ndarray and adds padding with zeros on right or bottom side of image (depends from aspect ratio), after that
|
||
converts data to float32 data type and change range of values from [0, 255] to [-1, 1], finally, converts data layout from planar NHWC to NCHW.
|
||
The function returns preprocessed input tensor and padding size, which can be used in postprocessing.
|
||
|
||
Parameters:
|
||
image (PIL.Image.Image): input image
|
||
Returns:
|
||
image (np.ndarray): preprocessed image tensor
|
||
meta (Dict): dictionary with preprocessing metadata info
|
||
"""
|
||
src_width, src_height = image.size
|
||
dst_width, dst_height = scale_fit_to_window(512, 512, src_width, src_height)
|
||
image = np.array(image.resize((dst_width, dst_height), resample=PIL.Image.Resampling.LANCZOS))[None, :]
|
||
pad_width = 512 - dst_width
|
||
pad_height = 512 - dst_height
|
||
pad = ((0, 0), (0, pad_height), (0, pad_width), (0, 0))
|
||
image = np.pad(image, pad, mode="constant")
|
||
image = image.astype(np.float32) / 255.0
|
||
image = 2.0 * image - 1.0
|
||
image = image.transpose(0, 3, 1, 2)
|
||
return image, {"padding": pad, "src_width": src_width, "src_height": src_height}
|
||
|
||
|
||
class OVStableDiffusionPipeline(DiffusionPipeline):
|
||
def __init__(
|
||
self,
|
||
vae_decoder: ov.Model,
|
||
text_encoder: ov.Model,
|
||
tokenizer: CLIPTokenizer,
|
||
unet: ov.Model,
|
||
scheduler: Union[DDIMScheduler, PNDMScheduler, LMSDiscreteScheduler],
|
||
vae_encoder: ov.Model = None,
|
||
):
|
||
"""
|
||
Pipeline for text-to-image generation using Stable Diffusion.
|
||
Parameters:
|
||
vae_decoder (Model):
|
||
Variational Auto-Encoder (VAE) Model to decode images to and from latent representations.
|
||
text_encoder (Model):
|
||
Frozen text-encoder. Stable Diffusion uses the text portion of
|
||
[CLIP](https://huggingface.co/docs/transformers/model_doc/clip#transformers.CLIPTextModel), specifically
|
||
the clip-vit-large-patch14(https://huggingface.co/openai/clip-vit-large-patch14) variant.
|
||
tokenizer (CLIPTokenizer):
|
||
Tokenizer of class CLIPTokenizer(https://huggingface.co/docs/transformers/v4.21.0/en/model_doc/clip#transformers.CLIPTokenizer).
|
||
unet (Model): Conditional U-Net architecture to denoise the encoded image latents.
|
||
vae_encoder (Model):
|
||
Variational Auto-Encoder (VAE) Model to encode images to latent representation.
|
||
scheduler (SchedulerMixin):
|
||
A scheduler to be used in combination with unet to denoise the encoded image latents. Can be one of
|
||
DDIMScheduler, LMSDiscreteScheduler, or PNDMScheduler.
|
||
"""
|
||
super().__init__()
|
||
self.scheduler = scheduler
|
||
self.vae_decoder = vae_decoder
|
||
self.vae_encoder = vae_encoder
|
||
self.text_encoder = text_encoder
|
||
self.unet = unet
|
||
self.register_to_config(unet=unet)
|
||
self._text_encoder_output = text_encoder.output(0)
|
||
self._unet_output = unet.output(0)
|
||
self._vae_d_output = vae_decoder.output(0)
|
||
self._vae_e_output = vae_encoder.output(0) if vae_encoder is not None else None
|
||
self.height = self.unet.input(0).shape[2] * 8
|
||
self.width = self.unet.input(0).shape[3] * 8
|
||
self.tokenizer = tokenizer
|
||
|
||
def __call__(
|
||
self,
|
||
prompt: Union[str, List[str]],
|
||
image: PIL.Image.Image = None,
|
||
negative_prompt: Union[str, List[str]] = None,
|
||
num_inference_steps: Optional[int] = 50,
|
||
guidance_scale: Optional[float] = 7.5,
|
||
eta: Optional[float] = 0.0,
|
||
output_type: Optional[str] = "pil",
|
||
seed: Optional[int] = None,
|
||
strength: float = 1.0,
|
||
):
|
||
"""
|
||
Function invoked when calling the pipeline for generation.
|
||
Parameters:
|
||
prompt (str or List[str]):
|
||
The prompt or prompts to guide the image generation.
|
||
image (PIL.Image.Image, *optional*, None):
|
||
Intinal image for generation.
|
||
negative_prompt (str or List[str]):
|
||
The negative prompt or prompts to guide the image generation.
|
||
num_inference_steps (int, *optional*, defaults to 50):
|
||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||
expense of slower inference.
|
||
guidance_scale (float, *optional*, defaults to 7.5):
|
||
Guidance scale as defined in Classifier-Free Diffusion Guidance(https://arxiv.org/abs/2207.12598).
|
||
guidance_scale is defined as `w` of equation 2.
|
||
Higher guidance scale encourages to generate images that are closely linked to the text prompt,
|
||
usually at the expense of lower image quality.
|
||
eta (float, *optional*, defaults to 0.0):
|
||
Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to
|
||
[DDIMScheduler], will be ignored for others.
|
||
output_type (`str`, *optional*, defaults to "pil"):
|
||
The output format of the generate image. Choose between
|
||
[PIL](https://pillow.readthedocs.io/en/stable/): PIL.Image.Image or np.array.
|
||
seed (int, *optional*, None):
|
||
Seed for random generator state initialization.
|
||
strength (int, *optional*, 1.0):
|
||
strength between initial image and generated in Image-to-Image pipeline, do not used in Text-to-Image
|
||
Returns:
|
||
Dictionary with keys:
|
||
sample - the last generated image PIL.Image.Image or np.array
|
||
"""
|
||
if seed is not None:
|
||
np.random.seed(seed)
|
||
# here `guidance_scale` is defined analog to the guidance weight `w` of equation (2)
|
||
# of the Imagen paper: https://arxiv.org/pdf/2205.11487.pdf . `guidance_scale = 1`
|
||
# corresponds to doing no classifier free guidance.
|
||
do_classifier_free_guidance = guidance_scale > 1.0
|
||
# get prompt text embeddings
|
||
text_embeddings = self._encode_prompt(
|
||
prompt,
|
||
do_classifier_free_guidance=do_classifier_free_guidance,
|
||
negative_prompt=negative_prompt,
|
||
)
|
||
# set timesteps
|
||
accepts_offset = "offset" in set(inspect.signature(self.scheduler.set_timesteps).parameters.keys())
|
||
extra_set_kwargs = {}
|
||
if accepts_offset:
|
||
extra_set_kwargs["offset"] = 1
|
||
|
||
self.scheduler.set_timesteps(num_inference_steps, **extra_set_kwargs)
|
||
timesteps, num_inference_steps = self.get_timesteps(num_inference_steps, strength)
|
||
latent_timestep = timesteps[:1]
|
||
|
||
# get the initial random noise unless the user supplied it
|
||
latents, meta = self.prepare_latents(image, latent_timestep)
|
||
|
||
# prepare extra kwargs for the scheduler step, since not all schedulers have the same signature
|
||
# eta (η) is only used with the DDIMScheduler, it will be ignored for other schedulers.
|
||
# eta corresponds to η in DDIM paper: https://arxiv.org/abs/2010.02502
|
||
# and should be between [0, 1]
|
||
accepts_eta = "eta" in set(inspect.signature(self.scheduler.step).parameters.keys())
|
||
extra_step_kwargs = {}
|
||
if accepts_eta:
|
||
extra_step_kwargs["eta"] = eta
|
||
|
||
for t in self.progress_bar(timesteps):
|
||
# expand the latents if we are doing classifier free guidance
|
||
latent_model_input = np.concatenate([latents] * 2) if do_classifier_free_guidance else latents
|
||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||
|
||
# predict the noise residual
|
||
noise_pred = self.unet([latent_model_input, np.array(t, dtype=np.float32), text_embeddings])[self._unet_output]
|
||
# perform guidance
|
||
if do_classifier_free_guidance:
|
||
noise_pred_uncond, noise_pred_text = noise_pred[0], noise_pred[1]
|
||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||
|
||
# compute the previous noisy sample x_t -> x_t-1
|
||
latents = self.scheduler.step(torch.from_numpy(noise_pred), t, torch.from_numpy(latents), **extra_step_kwargs)["prev_sample"].numpy()
|
||
# scale and decode the image latents with vae
|
||
image = self.vae_decoder(latents * (1 / 0.18215))[self._vae_d_output]
|
||
|
||
image = self.postprocess_image(image, meta, output_type)
|
||
return {"sample": image}
|
||
|
||
def _encode_prompt(
|
||
self,
|
||
prompt: Union[str, List[str]],
|
||
num_images_per_prompt: int = 1,
|
||
do_classifier_free_guidance: bool = True,
|
||
negative_prompt: Union[str, List[str]] = None,
|
||
):
|
||
"""
|
||
Encodes the prompt into text encoder hidden states.
|
||
|
||
Parameters:
|
||
prompt (str or list(str)): prompt to be encoded
|
||
num_images_per_prompt (int): number of images that should be generated per prompt
|
||
do_classifier_free_guidance (bool): whether to use classifier free guidance or not
|
||
negative_prompt (str or list(str)): negative prompt to be encoded
|
||
Returns:
|
||
text_embeddings (np.ndarray): text encoder hidden states
|
||
"""
|
||
batch_size = len(prompt) if isinstance(prompt, list) else 1
|
||
|
||
# tokenize input prompts
|
||
text_inputs = self.tokenizer(
|
||
prompt,
|
||
padding="max_length",
|
||
max_length=self.tokenizer.model_max_length,
|
||
truncation=True,
|
||
return_tensors="np",
|
||
)
|
||
text_input_ids = text_inputs.input_ids
|
||
|
||
text_embeddings = self.text_encoder(text_input_ids)[self._text_encoder_output]
|
||
|
||
# duplicate text embeddings for each generation per prompt
|
||
if num_images_per_prompt != 1:
|
||
bs_embed, seq_len, _ = text_embeddings.shape
|
||
text_embeddings = np.tile(text_embeddings, (1, num_images_per_prompt, 1))
|
||
text_embeddings = np.reshape(text_embeddings, (bs_embed * num_images_per_prompt, seq_len, -1))
|
||
|
||
# get unconditional embeddings for classifier free guidance
|
||
if do_classifier_free_guidance:
|
||
uncond_tokens: List[str]
|
||
max_length = text_input_ids.shape[-1]
|
||
if negative_prompt is None:
|
||
uncond_tokens = [""] * batch_size
|
||
elif isinstance(negative_prompt, str):
|
||
uncond_tokens = [negative_prompt]
|
||
else:
|
||
uncond_tokens = negative_prompt
|
||
uncond_input = self.tokenizer(
|
||
uncond_tokens,
|
||
padding="max_length",
|
||
max_length=max_length,
|
||
truncation=True,
|
||
return_tensors="np",
|
||
)
|
||
|
||
uncond_embeddings = self.text_encoder(uncond_input.input_ids)[self._text_encoder_output]
|
||
|
||
# duplicate unconditional embeddings for each generation per prompt, using mps friendly method
|
||
seq_len = uncond_embeddings.shape[1]
|
||
uncond_embeddings = np.tile(uncond_embeddings, (1, num_images_per_prompt, 1))
|
||
uncond_embeddings = np.reshape(uncond_embeddings, (batch_size * num_images_per_prompt, seq_len, -1))
|
||
|
||
# For classifier free guidance, we need to do two forward passes.
|
||
# Here we concatenate the unconditional and text embeddings into a single batch
|
||
# to avoid doing two forward passes
|
||
text_embeddings = np.concatenate([uncond_embeddings, text_embeddings])
|
||
|
||
return text_embeddings
|
||
|
||
def prepare_latents(self, image: PIL.Image.Image = None, latent_timestep: torch.Tensor = None):
|
||
"""
|
||
Function for getting initial latents for starting generation
|
||
|
||
Parameters:
|
||
image (PIL.Image.Image, *optional*, None):
|
||
Input image for generation, if not provided randon noise will be used as starting point
|
||
latent_timestep (torch.Tensor, *optional*, None):
|
||
Predicted by scheduler initial step for image generation, required for latent image mixing with nosie
|
||
Returns:
|
||
latents (np.ndarray):
|
||
Image encoded in latent space
|
||
"""
|
||
latents_shape = (1, 4, self.height // 8, self.width // 8)
|
||
noise = np.random.randn(*latents_shape).astype(np.float32)
|
||
if image is None:
|
||
# if we use LMSDiscreteScheduler, let's make sure latents are mulitplied by sigmas
|
||
if isinstance(self.scheduler, LMSDiscreteScheduler):
|
||
noise = noise * self.scheduler.sigmas[0].numpy()
|
||
return noise, {}
|
||
input_image, meta = preprocess(image)
|
||
latents = self.vae_encoder(input_image)[self._vae_e_output]
|
||
latents = latents * 0.18215
|
||
latents = self.scheduler.add_noise(torch.from_numpy(latents), torch.from_numpy(noise), latent_timestep).numpy()
|
||
return latents, meta
|
||
|
||
def postprocess_image(self, image: np.ndarray, meta: Dict, output_type: str = "pil"):
|
||
"""
|
||
Postprocessing for decoded image. Takes generated image decoded by VAE decoder, unpad it to initila image size (if required),
|
||
normalize and convert to [0, 255] pixels range. Optionally, convertes it from np.ndarray to PIL.Image format
|
||
|
||
Parameters:
|
||
image (np.ndarray):
|
||
Generated image
|
||
meta (Dict):
|
||
Metadata obtained on latents preparing step, can be empty
|
||
output_type (str, *optional*, pil):
|
||
Output format for result, can be pil or numpy
|
||
Returns:
|
||
image (List of np.ndarray or PIL.Image.Image):
|
||
Postprocessed images
|
||
"""
|
||
if "padding" in meta:
|
||
pad = meta["padding"]
|
||
(_, end_h), (_, end_w) = pad[1:3]
|
||
h, w = image.shape[2:]
|
||
unpad_h = h - end_h
|
||
unpad_w = w - end_w
|
||
image = image[:, :, :unpad_h, :unpad_w]
|
||
image = np.clip(image / 2 + 0.5, 0, 1)
|
||
image = np.transpose(image, (0, 2, 3, 1))
|
||
# 9. Convert to PIL
|
||
if output_type == "pil":
|
||
image = self.numpy_to_pil(image)
|
||
if "src_height" in meta:
|
||
orig_height, orig_width = meta["src_height"], meta["src_width"]
|
||
image = [img.resize((orig_width, orig_height), PIL.Image.Resampling.LANCZOS) for img in image]
|
||
else:
|
||
if "src_height" in meta:
|
||
orig_height, orig_width = meta["src_height"], meta["src_width"]
|
||
image = [cv2.resize(img, (orig_width, orig_width)) for img in image]
|
||
return image
|
||
|
||
def get_timesteps(self, num_inference_steps: int, strength: float):
|
||
"""
|
||
Helper function for getting scheduler timesteps for generation
|
||
In case of image-to-image generation, it updates number of steps according to strength
|
||
|
||
Parameters:
|
||
num_inference_steps (int):
|
||
number of inference steps for generation
|
||
strength (float):
|
||
value between 0.0 and 1.0, that controls the amount of noise that is added to the input image.
|
||
Values that approach 1.0 allow for lots of variations but will also produce images that are not semantically consistent with the input.
|
||
"""
|
||
# get the original timestep using init_timestep
|
||
init_timestep = min(int(num_inference_steps * strength), num_inference_steps)
|
||
|
||
t_start = max(num_inference_steps - init_timestep, 0)
|
||
timesteps = self.scheduler.timesteps[t_start:]
|
||
|
||
return timesteps, num_inference_steps - t_start
|
||
|
||
Configure Inference Pipeline
|
||
~~~~~~~~~~~~~~~~~~~~~~~~~~~~
|
||
|
||
|
||
|
||
First, you should create instances of OpenVINO Model.
|
||
|
||
.. code:: ipython3
|
||
|
||
import ipywidgets as widgets
|
||
|
||
core = ov.Core()
|
||
device = widgets.Dropdown(
|
||
options=core.available_devices + ["AUTO"],
|
||
value="AUTO",
|
||
description="Device:",
|
||
disabled=False,
|
||
)
|
||
|
||
device
|
||
|
||
|
||
|
||
|
||
.. parsed-literal::
|
||
|
||
Dropdown(description='Device:', index=4, options=('CPU', 'GPU.0', 'GPU.1', 'GPU.2', 'AUTO'), value='AUTO')
|
||
|
||
|
||
|
||
.. code:: ipython3
|
||
|
||
ov_config = {"INFERENCE_PRECISION_HINT": "f32"} if device.value != "CPU" else {}
|
||
|
||
text_enc = core.compile_model(TEXT_ENCODER_OV_PATH, device.value)
|
||
unet_model = core.compile_model(UNET_OV_PATH, device.value)
|
||
vae_decoder = core.compile_model(VAE_DECODER_OV_PATH, device.value, ov_config)
|
||
vae_encoder = core.compile_model(VAE_ENCODER_OV_PATH, device.value, ov_config)
|
||
|
||
Model tokenizer and scheduler are also important parts of the pipeline.
|
||
Let us define them and put all components together.
|
||
|
||
.. code:: ipython3
|
||
|
||
from transformers import CLIPTokenizer
|
||
|
||
scheduler = DDIMScheduler.from_config(conf) # DDIMScheduler is used because UNet quantization produces better results with it
|
||
tokenizer = CLIPTokenizer.from_pretrained("openai/clip-vit-large-patch14")
|
||
|
||
ov_pipe = OVStableDiffusionPipeline(
|
||
tokenizer=tokenizer,
|
||
text_encoder=text_enc,
|
||
unet=unet_model,
|
||
vae_encoder=vae_encoder,
|
||
vae_decoder=vae_decoder,
|
||
scheduler=scheduler,
|
||
)
|
||
|
||
Quantization
|
||
------------
|
||
|
||
|
||
|
||
`NNCF <https://github.com/openvinotoolkit/nncf/>`__ enables
|
||
post-training quantization by adding quantization layers into model
|
||
graph and then using a subset of the training dataset to initialize the
|
||
parameters of these additional quantization layers. Quantized operations
|
||
are executed in ``INT8`` instead of ``FP32``/``FP16`` making model
|
||
inference faster.
|
||
|
||
According to ``Stable Diffusion v2`` structure, the UNet model takes up
|
||
significant portion of the overall pipeline execution time. Now we will
|
||
show you how to optimize the UNet part using
|
||
`NNCF <https://github.com/openvinotoolkit/nncf/>`__ to reduce
|
||
computation cost and speed up the pipeline. Quantizing the rest of the
|
||
pipeline does not significantly improve inference performance but can
|
||
lead to a substantial degradation of accuracy.
|
||
|
||
For this model we apply quantization in hybrid mode which means that we
|
||
quantize: (1) weights of MatMul and Embedding layers and (2) activations
|
||
of other layers. The steps are the following:
|
||
|
||
1. Create a calibration dataset for quantization.
|
||
2. Collect operations with weights.
|
||
3. Run ``nncf.compress_model()`` to compress only the model weights.
|
||
4. Run ``nncf.quantize()`` on the compressed model with weighted
|
||
operations ignored by providing ``ignored_scope`` parameter.
|
||
5. Save the ``INT8`` model using ``openvino.save_model()`` function.
|
||
|
||
Please select below whether you would like to run quantization to
|
||
improve model inference speed.
|
||
|
||
**NOTE**: Quantization is time and memory consuming operation.
|
||
Running quantization code below may take some time.
|
||
|
||
.. code:: ipython3
|
||
|
||
to_quantize = widgets.Checkbox(
|
||
value=True,
|
||
description="Quantization",
|
||
disabled=False,
|
||
)
|
||
|
||
to_quantize
|
||
|
||
|
||
|
||
|
||
.. parsed-literal::
|
||
|
||
Checkbox(value=True, description='Quantization')
|
||
|
||
|
||
|
||
.. code:: ipython3
|
||
|
||
# Fetch `skip_kernel_extension` module
|
||
import requests
|
||
|
||
r = requests.get(
|
||
url="https://raw.githubusercontent.com/openvinotoolkit/openvino_notebooks/latest/utils/skip_kernel_extension.py",
|
||
)
|
||
open("skip_kernel_extension.py", "w").write(r.text)
|
||
|
||
int8_ov_pipe = None
|
||
|
||
%load_ext skip_kernel_extension
|
||
|
||
Prepare calibration dataset
|
||
~~~~~~~~~~~~~~~~~~~~~~~~~~~
|
||
|
||
|
||
|
||
We use a portion of
|
||
`conceptual_captions <https://huggingface.co/datasets/conceptual_captions>`__
|
||
dataset from Hugging Face as calibration data. To collect intermediate
|
||
model inputs for calibration we should customize ``CompiledModel``.
|
||
|
||
.. code:: ipython3
|
||
|
||
%%skip not $to_quantize.value
|
||
|
||
import datasets
|
||
import numpy as np
|
||
from tqdm.notebook import tqdm
|
||
from typing import Any, Dict, List
|
||
|
||
|
||
def disable_progress_bar(pipeline, disable=True):
|
||
if not hasattr(pipeline, "_progress_bar_config"):
|
||
pipeline._progress_bar_config = {'disable': disable}
|
||
else:
|
||
pipeline._progress_bar_config['disable'] = disable
|
||
|
||
|
||
class CompiledModelDecorator(ov.CompiledModel):
|
||
def __init__(self, compiled_model: ov.CompiledModel, data_cache: List[Any] = None, keep_prob: float = 0.5):
|
||
super().__init__(compiled_model)
|
||
self.data_cache = data_cache if data_cache is not None else []
|
||
self.keep_prob = keep_prob
|
||
|
||
def __call__(self, *args, **kwargs):
|
||
if np.random.rand() <= self.keep_prob:
|
||
self.data_cache.append(*args)
|
||
return super().__call__(*args, **kwargs)
|
||
|
||
|
||
def collect_calibration_data(ov_pipe, calibration_dataset_size: int, num_inference_steps: int) -> List[Dict]:
|
||
original_unet = ov_pipe.unet
|
||
calibration_data = []
|
||
ov_pipe.unet = CompiledModelDecorator(original_unet, calibration_data, keep_prob=0.7)
|
||
disable_progress_bar(ov_pipe)
|
||
|
||
dataset = datasets.load_dataset("conceptual_captions", split="train").shuffle(seed=42)
|
||
|
||
# Run inference for data collection
|
||
pbar = tqdm(total=calibration_dataset_size)
|
||
for batch in dataset:
|
||
prompt = batch["caption"]
|
||
if len(prompt) > ov_pipe.tokenizer.model_max_length:
|
||
continue
|
||
ov_pipe(prompt, num_inference_steps=num_inference_steps, seed=1)
|
||
pbar.update(len(calibration_data) - pbar.n)
|
||
if pbar.n >= calibration_dataset_size:
|
||
break
|
||
|
||
disable_progress_bar(ov_pipe, disable=False)
|
||
ov_pipe.unet = original_unet
|
||
return calibration_data
|
||
|
||
Run Hybrid Model Quantization
|
||
~~~~~~~~~~~~~~~~~~~~~~~~~~~~~
|
||
|
||
|
||
|
||
.. code:: ipython3
|
||
|
||
%%skip not $to_quantize.value
|
||
|
||
from collections import deque
|
||
from transformers import set_seed
|
||
import nncf
|
||
|
||
def get_operation_const_op(operation, const_port_id: int):
|
||
node = operation.input_value(const_port_id).get_node()
|
||
queue = deque([node])
|
||
constant_node = None
|
||
allowed_propagation_types_list = ["Convert", "FakeQuantize", "Reshape"]
|
||
|
||
while len(queue) != 0:
|
||
curr_node = queue.popleft()
|
||
if curr_node.get_type_name() == "Constant":
|
||
constant_node = curr_node
|
||
break
|
||
if len(curr_node.inputs()) == 0:
|
||
break
|
||
if curr_node.get_type_name() in allowed_propagation_types_list:
|
||
queue.append(curr_node.input_value(0).get_node())
|
||
|
||
return constant_node
|
||
|
||
|
||
def is_embedding(node) -> bool:
|
||
allowed_types_list = ["f16", "f32", "f64"]
|
||
const_port_id = 0
|
||
input_tensor = node.input_value(const_port_id)
|
||
if input_tensor.get_element_type().get_type_name() in allowed_types_list:
|
||
const_node = get_operation_const_op(node, const_port_id)
|
||
if const_node is not None:
|
||
return True
|
||
|
||
return False
|
||
|
||
|
||
def collect_ops_with_weights(model):
|
||
ops_with_weights = []
|
||
for op in model.get_ops():
|
||
if op.get_type_name() == "MatMul":
|
||
constant_node_0 = get_operation_const_op(op, const_port_id=0)
|
||
constant_node_1 = get_operation_const_op(op, const_port_id=1)
|
||
if constant_node_0 or constant_node_1:
|
||
ops_with_weights.append(op.get_friendly_name())
|
||
if op.get_type_name() == "Gather" and is_embedding(op):
|
||
ops_with_weights.append(op.get_friendly_name())
|
||
|
||
return ops_with_weights
|
||
|
||
UNET_INT8_OV_PATH = sd2_1_model_dir / 'unet_optimized.xml'
|
||
if not UNET_INT8_OV_PATH.exists():
|
||
calibration_dataset_size = 300
|
||
set_seed(1)
|
||
unet_calibration_data = collect_calibration_data(ov_pipe,
|
||
calibration_dataset_size=calibration_dataset_size,
|
||
num_inference_steps=50)
|
||
|
||
unet = core.read_model(UNET_OV_PATH)
|
||
|
||
# Collect operations which weights will be compressed
|
||
unet_ignored_scope = collect_ops_with_weights(unet)
|
||
|
||
# Compress model weights
|
||
compressed_unet = nncf.compress_weights(unet, ignored_scope=nncf.IgnoredScope(types=['Convolution']))
|
||
|
||
# Quantize both weights and activations of Convolution layers
|
||
quantized_unet = nncf.quantize(
|
||
model=compressed_unet,
|
||
calibration_dataset=nncf.Dataset(unet_calibration_data),
|
||
subset_size=calibration_dataset_size,
|
||
model_type=nncf.ModelType.TRANSFORMER,
|
||
ignored_scope=nncf.IgnoredScope(names=unet_ignored_scope),
|
||
advanced_parameters=nncf.AdvancedQuantizationParameters(smooth_quant_alpha=-1)
|
||
)
|
||
|
||
ov.save_model(quantized_unet, UNET_INT8_OV_PATH)
|
||
|
||
|
||
.. parsed-literal::
|
||
|
||
INFO:nncf:NNCF initialized successfully. Supported frameworks detected: torch, onnx, openvino
|
||
|
||
|
||
.. code:: ipython3
|
||
|
||
%%skip not $to_quantize.value
|
||
|
||
int8_unet_model = core.compile_model(UNET_INT8_OV_PATH, device.value)
|
||
int8_ov_pipe = OVStableDiffusionPipeline(
|
||
tokenizer=tokenizer,
|
||
text_encoder=text_enc,
|
||
unet=int8_unet_model,
|
||
vae_encoder=vae_encoder,
|
||
vae_decoder=vae_decoder,
|
||
scheduler=scheduler
|
||
)
|
||
|
||
Compare UNet file size
|
||
~~~~~~~~~~~~~~~~~~~~~~
|
||
|
||
.. code:: ipython3
|
||
|
||
%%skip not $to_quantize.value
|
||
|
||
fp16_ir_model_size = UNET_OV_PATH.with_suffix(".bin").stat().st_size / 1024
|
||
quantized_model_size = UNET_INT8_OV_PATH.with_suffix(".bin").stat().st_size / 1024
|
||
|
||
print(f"FP16 model size: {fp16_ir_model_size:.2f} KB")
|
||
print(f"INT8 model size: {quantized_model_size:.2f} KB")
|
||
print(f"Model compression rate: {fp16_ir_model_size / quantized_model_size:.3f}")
|
||
|
||
|
||
.. parsed-literal::
|
||
|
||
FP16 model size: 1691232.51 KB
|
||
INT8 model size: 846918.58 KB
|
||
Model compression rate: 1.997
|
||
|
||
|
||
Compare inference time of the FP16 and INT8 pipelines
|
||
~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~
|
||
|
||
|
||
|
||
To measure the inference performance of the ``FP16`` and ``INT8``
|
||
pipelines, we use median inference time on calibration subset.
|
||
|
||
**NOTE**: For the most accurate performance estimation, it is
|
||
recommended to run ``benchmark_app`` in a terminal/command prompt
|
||
after closing other applications.
|
||
|
||
.. code:: ipython3
|
||
|
||
%%skip not $to_quantize.value
|
||
|
||
import time
|
||
|
||
def calculate_inference_time(pipeline, validation_data):
|
||
inference_time = []
|
||
pipeline.set_progress_bar_config(disable=True)
|
||
for prompt in validation_data:
|
||
start = time.perf_counter()
|
||
_ = pipeline(prompt, num_inference_steps=10, seed=0)
|
||
end = time.perf_counter()
|
||
delta = end - start
|
||
inference_time.append(delta)
|
||
return np.median(inference_time)
|
||
|
||
.. code:: ipython3
|
||
|
||
%%skip not $to_quantize.value
|
||
|
||
validation_size = 10
|
||
validation_dataset = datasets.load_dataset("conceptual_captions", split="train", streaming=True).take(validation_size)
|
||
validation_data = [batch["caption"] for batch in validation_dataset]
|
||
|
||
fp_latency = calculate_inference_time(ov_pipe, validation_data)
|
||
int8_latency = calculate_inference_time(int8_ov_pipe, validation_data)
|
||
print(f"Performance speed-up: {fp_latency / int8_latency:.3f}")
|
||
|
||
|
||
.. parsed-literal::
|
||
|
||
/home/nsavel/venvs/ov_notebooks_tmp/lib/python3.8/site-packages/datasets/load.py:1429: FutureWarning: The repository for conceptual_captions contains custom code which must be executed to correctly load the dataset. You can inspect the repository content at https://hf.co/datasets/conceptual_captions
|
||
You can avoid this message in future by passing the argument `trust_remote_code=True`.
|
||
Passing `trust_remote_code=True` will be mandatory to load this dataset from the next major release of `datasets`.
|
||
warnings.warn(
|
||
|
||
|
||
.. parsed-literal::
|
||
|
||
Performance speed-up: 1.232
|
||
|
||
|
||
Run Text-to-Image generation
|
||
----------------------------
|
||
|
||
|
||
|
||
Now, you can define a text prompts for image generation and run
|
||
inference pipeline. Optionally, you can also change the random generator
|
||
seed for latent state initialization and number of steps.
|
||
|
||
**Note**: Consider increasing ``steps`` to get more precise results.
|
||
A suggested value is ``50``, but it will take longer time to process.
|
||
|
||
Please select below whether you would like to use the quantized model to
|
||
launch the interactive demo.
|
||
|
||
.. code:: ipython3
|
||
|
||
quantized_model_present = int8_ov_pipe is not None
|
||
|
||
use_quantized_model = widgets.Checkbox(
|
||
value=True if quantized_model_present else False,
|
||
description="Use quantized model",
|
||
disabled=not quantized_model_present,
|
||
)
|
||
|
||
use_quantized_model
|
||
|
||
|
||
|
||
|
||
.. parsed-literal::
|
||
|
||
Checkbox(value=True, description='Use quantized model')
|
||
|
||
|
||
|
||
.. code:: ipython3
|
||
|
||
import gradio as gr
|
||
|
||
|
||
pipeline = int8_ov_pipe if use_quantized_model.value else ov_pipe
|
||
|
||
|
||
def generate(prompt, negative_prompt, seed, num_steps, _=gr.Progress(track_tqdm=True)):
|
||
result = pipeline(
|
||
prompt,
|
||
negative_prompt=negative_prompt,
|
||
num_inference_steps=num_steps,
|
||
seed=seed,
|
||
)
|
||
return result["sample"][0]
|
||
|
||
|
||
gr.close_all()
|
||
demo = gr.Interface(
|
||
generate,
|
||
[
|
||
gr.Textbox(
|
||
"valley in the Alps at sunset, epic vista, beautiful landscape, 4k, 8k",
|
||
label="Prompt",
|
||
),
|
||
gr.Textbox(
|
||
"frames, borderline, text, charachter, duplicate, error, out of frame, watermark, low quality, ugly, deformed, blur",
|
||
label="Negative prompt",
|
||
),
|
||
gr.Slider(value=42, label="Seed", maximum=10000000),
|
||
gr.Slider(value=25, label="Steps", minimum=1, maximum=50),
|
||
],
|
||
"image",
|
||
)
|
||
|
||
try:
|
||
demo.queue().launch()
|
||
except Exception:
|
||
demo.queue().launch(share=True)
|